# Stable Diffusion Optimization

This fine-tuned model optimizes short prompts into detailed ones, enabling image generation based on Stable Diffusion or Flux technology.

A text-to-image generation model based on Stable Diffusion XL (SDXL), supporting ControlNet functionality.

A text-to-image generation model based on Stable Diffusion XL, supporting extended features like ControlNet.

A text-to-image generation model based on Stable Diffusion XL, supporting ControlNet functionality.

Anti-blur image generation model based on Stable Diffusion and LoRA technology

LoRA model trained on 25 photos of Livvy Dunne for text-to-image generation tasks

This is a LoRA model based on Stable Diffusion technology, specifically designed for generating images in the style of Elle Fanning.

Trained on 25 portrait photos of Gal Gadot, using SimpleTuner tool for 1600 iterations of LoRA model

This is a LoRA model based on Stable Diffusion technology, specializing in generating Pixar-style 3D cartoon images.

This is a text-to-image generation model combining Hands XL, SD 1.5, and FLUX.1-dev technologies, focusing on high-quality image generation.

A 3D animation-style text-to-image generation model based on Stable Diffusion technology, supporting Disney and anime-style 3D image generation.

A LoRA model trained on 25 images of Jenna Ortega, using SimpleTuner for 1600 steps, capable of generating relevant images without trigger words.

LoRA model trained on 25 photos of Emma Watson for generating Emma Watson-style images

A text-to-image generation model fine-tuned based on Stable Diffusion v1.5, using aesthetic reward functions on animal datasets and trained with reward backpropagation methods.

Lexica Aperture V3.5 is a free text-to-image generation model, similar to the paid version of Lexica Aperture V3.5, supporting multilingual input.

A LoRA model for sticker image generation compatible with SD 1.5, capable of creating anime sticker images in various styles

A set of creative embedding models for Stable Diffusion designed to enhance image generation quality and support specific styles and effects.

A collection of anime-style LoRA fine-tuning models created by Hotaru Jujo, focusing on facial expressions and special effects in text-to-image tasks

A ControlNet model trained on Stable Diffusion v1.5, utilizing a novel conditional control method for image-to-image conversion

A photorealistic model based on Stable Diffusion 1.5, focused on generating high-quality, realistic images.

Anything V3 is a text-to-image generation model based on Stable Diffusion technology, capable of generating high-quality images from textual descriptions.

This is a fashion model image generation model trained using Dreambooth technology, capable of generating high-quality fashion model images based on text prompts.

A text-to-image generation model trained using TheLastBen's fast-DreamBooth notebook, supporting Stable Diffusion technology.

ProllyMix is a text-to-image generation model based on Stable Diffusion, suitable for creative image generation tasks.

A realism art-style image generation model fine-tuned based on Stable Diffusion 1.5

A text-to-image model based on Stable Diffusion technology, capable of generating high-quality images from text descriptions

A text-to-image model trained by michaelkmccoy using buildspace's DreamBooth notebook, suitable for creating personalized AI avatars

A diffusion model for generating modern art style images

FurryDiffusion is a Stable Diffusion model specifically designed for generating furry art, currently in beta testing and continuously being optimized.

A Dreambooth fine-tuned model based on Stable Diffusion for text-to-image generation tasks

This is a text-to-image generation model based on Stable Diffusion technology, trained using the DreamBooth method, specifically designed for generating sci-fi style landscape images.

A Stable Diffusion model trained with Dreambooth technology for generating images of specific styles

A Stable Diffusion model trained with Dreambooth technology for generating images in specific styles

An anime-style pencil concept sketch generation model fine-tuned based on Stable Diffusion 1.5

Text-to-image generation model based on Stable Diffusion technology, trained using the Dreambooth method

An AI model for generating cyberpunk-style anime characters, developed based on fine-tuned Waifu Diffusion V1.3 and Stable Diffusion V1.5

A concept model fine-tuned using Dreambooth technology based on the Stable Diffusion 1.4 framework for the movie version of Tails from 'Sonic the Hedgehog 2'